# Creative Image Generation

This is a LoRA adapter trained for the FLUX.1-dev text-to-image model, supporting invocation via diffusers or ComfyUI to generate images of specific styles.

FLUX.1-dev is a text-to-image model based on double exposure technology, specializing in the artistic fusion of human silhouettes with various scenes.

This is a text-to-image generation model based on a diffuser, capable of generating corresponding images from input text descriptions.

LoRA fine-tuned model based on Stable Diffusion XL, specialized in text-to-image tasks

A LoRA fine-tuned model based on Stable Diffusion XL, specifically designed for generating ASCII art-style images.

DreamShaper is a powerful generative model focused on image generation and editing tasks.

ControlNet v1.1 is a neural network structure developed by Lvmin Zhang, designed to control large pre-trained diffusion models through additional conditions. This version is trained based on instruct pix2pix image conditions.